#💬|general-chat
1 messages · Page 77 of 1
you've got nvidiia so it'll work. i was just saying it would just not load if you tried in the case of getting an amd card.
if you're going from nvidia to nvidia, i'd just refresh drivers and see if it works
oh okay. well thank you
2 monrha with SDXL, still no Controlnet. What gives?
thank you for your reply. I have similar thoughts
HI guys, what do you think of Stable Audio. As an audio producer it's really impressing me. I hope the audio quality improves soon, the tracks all sound distorted. It would be good to hear your thoughts.
hi guys, i have a problem with the controlnet in the img2img tab, i canno't preview the preview, what can i do?
Stable Audio techies, we need the music output to be "flat" i.e. flat EQ, and check that output volume, bring it down a bit.
I was ruining my images. Clipdrop downsamples when you use the eraser. The site upscale.media upscales with a higher quality than clipdrop for free
you never know what is free. You can upscale on your owns in super quality and proper models. Using some this way oriented software
#✨|sdxl message
Check this @fervent thunder
hi .. using leonardo,.... I want to create two characters .. one with short hair and his girlfriend with long hair, now my prompt is this : "show two characters walking next to each other, max the adventurous boy with short dark brown hair, mia his girlfriend, orange red long hair" I am wasting really a lot of token generating two the same looking characters .. both have same hair lenght .. their faces are often the same and its ignoring my description .... could you tell me please what I do wrong ? ( I tried more descriptive prompts but nothing works ) ... help :))
maybe 1boy with short hair and 1girl with long hair. But you need probably perform it many times. If you were using A1111 i would suggest you regional prompter, where you divide screen, and then tell what is in what part. @slate heron
thanks @plucky shoal
for nothing 🙂 hope it will work. I know it works for some those 1boy and 1girl, but bit of luck is needed
ugh that worked at once ! ... GREAT !! thanks again @plucky shoal
np happy it helps
Yo
need to know something, is it possible to create a video (movie /reel kind of using deforum sd? auto111?
suppose i have a story "turtle jumps from parachute" or "naruto escaping from jail" (in a cinematic style?)
@plucky shoal hi, i remember you've helped me a lot with sd, is it possible^
Hi i am not sure, i only did one prompt to another prompt to another prompt. Like zoom in bit, or like fly througth.
better to ask somebody with video exp. There is site that doing very short videos but with say wink eye, so i think it is not still possible.
Will search for that site, somehow pico or so...
#✨|sdxl message
pikalabs
stv qlq peut m'aider ?
je doit trouver le "HIKS" en Gras double souligné italique, en cyan, avec la police Time New Roman, taille 4 STP, mais je galère un peut, voila le lien :
https://docs.google.com/spreadsheets/d/1aSLLN-1WrwuUB_dacDiuxh1QQ9HzUBBbkOra8kwmrUY/edit#gid=0
whats the prompt for the ai to use literal text?
as in i want "ice cream" to be written out in full in the image
stable diffusion really sucks at doing text. sdxl has some capabilities, but to prompt specifically for written text you probably thinking of deep floyd
isnt there a bracket mix that induces the gen to write out the text in full?
in deep floyd yes
whats that?
ok
https://huggingface.co/stabilityai/stable-diffusion-xl-base-1.0#limitations here's the sdxl model card. you'll note that in the limitations area, it mentions text.
Project address: https://github.com/CloudOrc/SolidUI
hello everyone,I'm working full-time on the open-source project SolidUI, which can generate any graph with just one sentence. I welcome friends to participate and contribute together.

Stable Stage is Live! 🔴
Hello. I consistently get BIAS words in negative embedding. This is BAD. Is there a way I can learn how to generate my own negative embedding?
What is Center Stage?
Anyone know what Würstchen is and how to use it?
The huggingface page has this: "Otherwise, the model is available through the Diffusers Library, so you can use the interface you are already familiar with. For example, this is how to run inference using the AutoPipeline:"
what in the world does that mean?
Ok i need some help
Can you change the perspective of a character's face if you use img2img?
I have a very weird request... I am attempting to use image to image to turn a... strange... photo into a renaissance painting that maintains the composition and detail of the original image and just makes it renaissance style... so far been unable to achieve that... Any tips for checkpoints or Lora that could help me out here?
The image in question is a uh... banana with a face
making it as a joke for my friend. And yes I mean an actual banana.. well what looks like a porceline one or something but
Tried uisng img2img with no luck,,, keeps turning the banan into a human
#🏞|general-with-images message image posted here
Instead of what im after im getting some nightmare fuel 
can someone help me with instalation ? i got to Step 5: Setup the Web-UI on this tutorial - https://www.datacamp.com/tutorial/how-to-run-stable-diffusion
Is there another way to use stable diffusion on your cell phone without collab?
I'm trying in all colabs but I'm not succeeding, I enter but 5 minutes later it gives an error
hi guys, I'm using Local stable diffusion installation and I have a problem with generating images because I wanted to generate image in 2560x2560 resolution but the highest resolution I can go for is 1024x1024 for img2img, every value over that gives me an error:
OutOfMemoryError: CUDA out of memory. See documentation for Memory Management and PYTORCH_CUDA_ALLOC_CONF.
I have INNO3D GeForce RTX 3060 Twin X2 LHR 12GB
What local installation of Stable Diffusion should I use if i wanna generate 2560x2560 resolution images and why img2img version has maximum 2048x2048 resolution. And what is the best way to get highres images without upscaling it. Or maybe the is good upscaling option I don't know about, because I was trying to upscale my 1024x1024 images but it's not that good as highres generated images.
As I'm a beginner to SD I'm a little bit confused... I'v installed the A1111 and some different models and LoRAs from civitai. But even if I use the same parameters, seed and so one from the example pictures my portraits look like mutants... need some help here 😦
can you give example pictures
Ok... - I've downloaded the Anya Taylor-Joy LoRA from civitai and also the models, used the same prompt and seed as on the site and I get this
post in the general-with-images channel
done
I've been messing around with Stable Diffusion for a few days and I can't seem to get my images to not look like crap. Body parts everywhere, art style looks like a kindergartner drew it even if I copy detailed prompts from a site
Any chance someone can chat me while I do a screen share and walk me through what I might be doing wrong? Certain settings etc
Upscaling black-and-white line art with R-ESRGAN 4x+. The upscaling is fairly flawless. I haven't really played around with other upscalers, but have noticed that there are more options now than there were six months ago. Are there any that you prefer instead of R-ESRGAN 4x+ for upscaling black-and-white line art?
Have you been using some negative prompts, negative LORA's, especially with anything based on 1.5?
Where in easy diffusion can i select the strength for image to image, of how much it looks like the original pic?
I haven't, are the negative prompts typically included in the examples such as lexica.art? I also don't know what LORAs are or 1.5
Are you using Automatic1111 or another platform?
I downloaded the checkpoint v1-5-pruned-emaonly.ckpt that the guide recommended
A1111 yes
although I don't know how to verify that, I just followed the steps for an a1111 guide
I don't know how to tell, is A1111 the model?
A1111 is the graphic user interface for the general program stable diffusion. You can use standard models or you can use other models.
The checkpoints are models.
I'm not sure what that model is then, I guess a generic one? v1-5-pruned-emaonly.ckpt
is lama cleaner slow? I'm using accelerate and cuda. With SD1.5 model it takes so long to erase the moon from the sky. 2% 4% 6% ....
Yep. that's the base model. Lots of friendly folks have created adapted models that are... better.
They're available on civitai.com, along with LORA's which are like plugin "helper" models that are smaller.
There are so many models! Any recommendations?
My favs: Realistic Vision V5.1, Dreamlike Photoreal 2.0, Any of the Dreamshaper versions, Deliberate... There are a lot. Some of them require a yaml file too, so watch for that.
Best overall? Tie between Dreamshaper 8 and Realistic Vision 5.
Download the model to whereever your A1111 install is here:
D:<your file system>\automatic1111\stable-diffusion-webui\models\Stable-diffusion
and your LORA's here:
D:<your file system>\automatic1111\stable-diffusion-webui\models\Lora
stuff you see on lexica could be from 1.2 or 1.4
so the same prompt on 1.5 wouldn't give you the same result
Reboot your A1111 and the new models/loras will become available.
is gpu faster than cpu to erase using AI?
yes
took 7 minutes to erase the moon
Do you have a text file with common negative prompts? Or do the models have the negative prompts build in
Where in easy diffusion can i select the strength for image to image, of how much it looks like the original pic?
Typical negative prompt example:
(deformed, distorted, disfigured:1.3), poorly drawn, bad anatomy, wrong anatomy, extra limb, missing limb, floating limbs, (mutated hands and fingers:1.4), disconnected limbs, mutation, mutated, ugly, disgusting, blurry, amputation, crossed eyes, bad irises, bad pupils, bad eyebrows, umbrella, long body, too many arms, umbrellas:5, stretched body:5
to reduce the likelihood of text I also use this in the negative prompt:
Kanji, Hiragana, Katakana, alphabet, syllabet, Cuneiform, Heiroglyphs, Ivrit, Rashi, Chinese Characters, Hebrew Characters, Hieratic, Demotic
This is mostly useful for 1.5 based models. Less so for XL models.
Also for stopping those crummy watermarks and signatures:
watermark, signature, logo, chop
And so do you keep text handy in an editor to copy and paste each time, or are these included in models?
Definitely keep notes and save them. There are also styles within the interface under the big orange generate button. That's a whole thing you could explore because it's very useful for quickly adding positive and negative prompts. You can also set A1111 to add the prompts you've been using to your saved pngs. If you lose your prompt, just find a previous image, drag it to the "PNG Info" tab, click either "send to txt2img" or the like. It will populate all the field for you. Remember to click the "random" (dice icon) button so you generate new stuff rather than that seed.
So nobody knows how to set the strength of how much the image to image should look like the original, in easy diffusion?
Sorry... A lot of us are using Automatic1111. I haven't taken the time to download Easy. I hope someone comes along to help.
Is automatic still that annoying to install. Or can i skip most steps because i already installed easy diffusion?
2560x2560 straight generations on 12gb, not going to happen. you're going to want to do a low resolution pass, 768 perhaps, then scale that up using an upscale tile script
i dont know how easy diffusion works, it's acronym, ED, is unfortunate though. What i do know is how img2img works. it's not a matter of making an image look closer to the original or less. It's how much noise you're adding to the image before processing it. on other UI's you see a denoise setting. the higher you tell it to denoise an image, the more noise it adds into it at the start. So higher denoise setting means less of the original image will be there.
there's a sweet spot around 0.4 - 0.6 denoising where images retain a lot of their original qualitys while still allowing diffusion to do some magic
Using online stable diffusion models i could set in % how close it would look to the original
I did "git clone https://github.com/AUTOMATIC1111/stable-diffusion-webui.git" and it downloaded it...but where exactly did it download it to?
@pale latch so I need to generate in resolution for example 1024x1024 and then upscale my image with some extentions?
that'd be more reaosnable. i'm not sure where you got the idea that 12gb would be capable of shooting out 5mp images
2048x2048 is the near limit of my 16gb card
?
Run it on another remote machien/service and connect your phone to that interface
or dreambooth or civitai or a hundred other web service SD pages
I'm new in this, how do I do this?
starts with being willing to pay money
Isn't there a free way?
collab was the free service provided by google, but 1000s of horny teens abused it for porn and now the free lunch is over
100,000s more accurately
I did hentai XD
google probably would've kept it free if people were just using the service for learning like it's intended
but, people did not
i myself benefitted from collab a lot while i was leraning. i didn't use the UIs though. They weren't available on collab yet. I used scripts and just edited them to see what they did. Cool stuff. I ended up getting it working on my AMD card in linux and then buying hardware once i had figured things out enough
point is, once i had learned enough, i stopped using collab. that's what it's there for
all the people who just abused the shit out of the service instead of using it for learning purposes, are the reason google is banning accounts that use webui's now
it's quite a shame that all new users now can't enjoy that boost google was offering before
I think I'm a little to blame for this too, I used stable diffusion to create hentai, not just hentai but I created
https://youtu.be/MaSRXvM05x4 here's a cool service being offered by civit. hook their generation servers up to your own local installed A1111 webui. but youv'e gotta pay
"Pornography drives markets" is often heard a lot, but in this case it actually ruined all the free access people had before. Nobody wants to offer free generation services or even touch the idea with a 100 foot pole
I don't have any money, unfortunately
The short window where this was a free hobby is over. You could try looking into something like "stable horde", a group of volunteers providing generation services to people for free
but it's a token based system and they're treated like crypto currency that the people who run the network centrally control
so basically, a paid saas like any other service
Is there a way to use free stable diffusion on PC?
by hobby you mean playing around with AI art?
thats right
its still free if you use a gpu you own, but that's a cost of entry situation
oh, yes thats true
but thats where Midjourney comes into the game as well
and for another target groups, Adobe Firefly
Weird with webui stable diffusion i get much better and faster results than with easy stable diffusion, even thru i use the same models
Right. I'm not saying that all image generation is a hobbyist situation. People who are seeking free services are assumably hobbyists though, imo
How i get control net on the webui?
yeah id say hobbyists are the ones who are heavily into free services. I see this in 3D spectrum with Blender
although im not a pro yet but still pay for 3ds Max and other non-free programs
but usually hobbyists are running for free services...at first
on the other hand, a lot of the people if not the absolute majority on Midjourney are hobbyists as well
part of them tries to earn money ofc with MJ
i used to be a pro in the graphic design field. i left and started working in other fields. now this has happeend so i'm a hobbyist with pro experience now. i started on the free collab scripts, no ui. then i got automatic1111 working in linux with an amd 8gb card. that worked out so well and i learned so much that when it really boosted me when i bought a new 4080
understand and im always interested in how people got to their workflow and what their workflow is in the first place 😄
i started with Crayion/DALL-E Mini back then XD
i mean gen AI
is this a scam?
I’m interested in buying a few of your artworks for my exhibition, about 5-10 artworks. But I’m only interested in buying them as NFTs.
I’ll offer $2500 for each of these artworks I purchase from you as I’d love to purchase these artworks as digital artworks(NFTs).
2500$ for randomly generated images for free with SDXL sounds way too good to be true.
scam
out of the blue someone in deviant art wants to buy my AI generated artwork offering over 1000$ ? wtf?
How would i install a sdxl model? Do i just make a new folder and put it in?
Or do i need to do something to even have sd xl?
anyone here know if there are any large collection of Styles.cvs files freely available?
tell me a common fix for the anime hands in different angles
without using controlnet
Does anyone here know how to speed up the iterations per second on a AMD GPU (RX 6800XT), using The Automatic1111 SD setup?
What do people on reddit mean anything but 1.5 is censored?
I tried some XL models, and i can make anything with them?
they mean that some of the filtering on the dataset was focused on removing pornographic images so they throw their arms up and declare oppressive censorship. it's like, because simpsons won't do porno episodes, that's an evil level of censorship
those same people have also started to shut up a lot more and only the dimmer ones are still going on about xl being censored
not porn or some anime styles,the anime styles cannot be done yet because of another reason not related to censorship
You can do nude people tho
nudity is not porn
I see
sdxl is version 1.0 too. sd 1.5 is a few iterations of the base. People can use the same dataset that created sd 1.5 to refine sdxl
So what is even special about XL?
more parameters. higher attention resolution. 2 clip layers. lots
Attention resolution?
yes we can we just need to replace the NAI dataset that contained like 5million images with a dataset of the same size so we just need a bunch of H100 or A100's to finetune sdxl
easy money
sd 1.5 is trained with 512x512 resolution, so it has that much attention space.
xl is 1024x1024
But it won't make it easier for me to make pics bigger than 512?
novel ai didn't make 1.5. runway ml did. and it only needs to be done once if you want to do it. they used the laion aesthetic set. Literally open and available to anyone
never said nai made 1.5 i said nai leak that had like 5million or 17million images is what made 1.5 anime models capable of doing the variety of anime styles it can do
the pornography is censored loudmouths are so silly. y'all paid unstability.ai to make all these community porno models and they're squandering all the cash on their generation service they sell. Their obvious trademark infringement shows that they're unwilling to comply with etiquette and rules, yet they raise nealry half a million
NansException: A tensor with all NaNs was produced in Unet. This could be either because there's not enough precision to represent the picture, or because your video card does not support half type. Try setting the "Upcast cross attention layer to float32" option in Settings > Stable Diffusion or using the --no-half commandline argument to fix this. Use --disable-nan-check commandline argument to disable this check.
What that means?
either add --no-half or change vae
personally i think xl has better anime capabilities out of the gate than nai has
and novelai used a publical booru dataset too
RuntimeError: mat1 and mat2 shapes cannot be multiplied (616x2048 and 768x320)
yea the only difference is size of the dataset if we could finetune on a 3090 we would probably reach the same levels of quality but too expensive to finetune sdxl
Model refining on a 3090 is expensive for 1.5 too. That's why its done on clusters generally.
Anyone can do this. I bet someone legitimately doing a project with real research that sought to bring superior anime capability could petition stability to donate cluster time. A lot of community good will would be available to them if they proved they were doing the work.
When i use control net, what impact does denoising strength have?
nobody is going to give money without seeing a dataset being worked on at this point. Not after Unstability swindled those good will patrons.
The opportunity is there for anyone to explore. I'd even expect they take some of the donations for themselves to cover the expenses of their time and energy. Would build a good brand for themselves too
Sucks about unstabiliyt, but they at least proved just how much money is out there willing to donate to this kind of cause
Like if i put normal map in net control, and check pixel perfect, what use does that have if it won't really follow the normal map if i put denoising to 100%?
pixelperfect is just for the preprocessor's benefit. if you're putting a normal map directly into controlnet, don't select any preprocessor
How you use the recolor net control thing?
and do i have to select any model for it?
these are advanced topics. helps to follow guides.
maybe don't jump into all the contorlnet models immediatly. learn the limits of prompting
Where i find some guide or something that explains all the 14 control net models?
So for example i use lineart, it guesses the lineart from a reference picture? But does it also make sense to select lineart if i want to turn already existing lineart into a picture?
lineart is only for coloring existing lineart into a picture
But the article says "Line Art renders the outline of an image. It attempts to convert it to a simple drawing."
And when you use it, it does that?
canny works better for converting a real pic into lineart like this video https://www.youtube.com/watch?v=xlO1Av20-OY
Like i just used it, and it also gave me the lineart pic with the generated pics
But what if i want to turn lineart into a picture instead?
u can use either canny or lineart try with both
So they both spit out lineart guessed from a picture for me, but they also can turn lineart into a picture?
just put the lineart pic into controlnet and enable it and play with the weights to see which gives u better results
yea u can also use a 2nd picture with color to give it that color to the generated image
this is another method that inverts the lineart pic https://www.youtube.com/watch?v=kZtoBSDUdEk
you'd select lineart controlnet model, and none for preprocessor
any model can take a raw input that you've given it. preprocess just works the reference image into what the model knows
try the lineart model with coloring book pages
I noticed that when the bot says that I need to vote more, if I try another channel, it usually works
I want to take one picture of one person, and one picture of another person, and use img2img or something similar to put them in an image together. is this possible?
I got this QR code control network thing, but somehow the reference thing i put in barley has any effect on the output image
how can i stop sd from saving grids in the output folder?
Settings
as stable diffusion programs can generate exactly same batch again and again so I think it would be great to have an option to set like a "batch range".. I mean like if i'm generating 20 images in a batch now later I can continue from 21 to upwards [like 21 to 40] instead of starting the batch from 0 to 20 same ones again. it would save us lots of time and power to gather new generations.
Hey sorry to disturb, I am a SD nerd , do u know a way to change emotions in a existing image? I want it for a yt video thumbnail and dont know enough photoshop
hello. I am trying to invite friend in this server, but the link is not working
does anyone the name of the extension that adds a vae dropdown menu next to the checkpoint dropdown menu at the top of a1111 webui?
it's in the settings.. look for the quick settings list
it says the invite link is expired.
oohhhh thank you
any suggestions on what else to add to the quick list thats useful?
How do people get 20 it/s on the gtx 4070 ti? I only get around 1-5. What do i do wrong?
Weird earlier i got 10
Ok on base i get 14 it/s
Not sure if its worth it optimizing for the last 4-5, and then might get worse results
maybe resolution, what UI are you using? ComfyUI or A1111?
can anyone explain how can you train AI with other pictures (i got a folder ready and installed SD) but i do not know how to use it
I don't even remember which one i installed
Oh automatic
I thought i got comfy
Does sdxl not work with image to image yet?
Same with netcontrol?
anybody using this? https://huggingface.co/stabilityai/stable-diffusion-x4-upscaler
@fervent thunder Have you --xformers. in your webui-user.bat?
Als speed depends on used sampler
I don't. I just have @echo off
set PYTHON=
set GIT=
set VENV_DIR=
set COMMANDLINE_ARGS=
call webui.bat
add --xformers in set COMMANDLINE_ARGS=
set COMMANDLINE_ARGS= --xformers
What will that do?
Do i have to edit the shell webui user file too?
Also i just delete set comandline args and replace it with
"--xformers in set COMMANDLINE_ARGS=
set COMMANDLINE_ARGS= --xformers" ?
Or do i just ad xformers after the equal symbol?
yes
Do i have to start it thru the bat file then? Or is the bat file just the settings for it?
Oh its installing xformers
Thanks
just restart if you are in. Run ...
Seems o.k. then, report what gain you have 🙂
🙂
Does stability ai have a public roadmap for the next few years?
hey, on DreamStudio I get the message: Something isn't quite right with your prompts.
with this prompt:
A bright-eyed, innocent young girl with chestnut hair in two pigtails, holding tightly onto a tattered plushy bear as she navigates the impossibly large and cave. The mood is mysterious, but her blooming resolve shines stronger than her fear. Illuminated by the faint glow of bio-luminescent mushrooms, the scene presents a narrative of bravery.
negative prompt: modern urban background ,adult characters, darker tones ,gore, warfare elements
prompt strength: 1/auto
generation steps:25
model: SDLX v0.9 and 1.0.
Can someone more knowledgable tell me why ?
Is it normal that once you turn hires fix on you get sometimes a entirly different pic, even with same seed?
What is the solution, have the images right from the start be done with hires, to get the end results already?
yes it is normal. Solution is using resizers from extra tab.
Is it possible to create SDXL textual inversion embeddings?
Would it be okay if someone could extend an image I have to fit into a phone wallpaper
How do I upload a custom image and generative fill it in the desired directions?
Can you already integrate animategif in a1111?
You need to lower the denois for it to look like the image without hires. Like you said you can just generate with highres fix on at start
I want to take one picture of one person, and one picture of another person, and use img2img or something similar to put them in an image together. is this possible?
RuntimeError: Expected weight to be a vector of size equal to the number of channels in input, but got weight of shape [1280] and input of shape [2, 2560, 8, 8]
Whats that? I put a image in animated diff, and i get the error, and after that i get the error every time now
I get that bug where on second run of trying to make a gif, it does that
Im not sure what to do. Did disable xformers, and the models are 100% in the right folder
On animatediff
Hello everyone, I need to make a cover for my Facebook, what measurements should I use in the image prompt?
for a profile pic square base res 512x512 and then upscale to 1024x1024
Thank you so much🥰
HellO!
Hi, I have a question. I downloaded one of the models and I want to use it with ComfyUI, and I have a pretty decent PC. My question is, is there a chance it will damage my PC/GPU? from overheating or something?
just get HWINFO and if your hardware reaches high temps it throttles anyway
New shared API key: https://tinybots.net/artbot?shared_key=f9821459-e5f6-4864-a8c3-749d936fb75d
The AI Horde section in #1080946152318443610 links to ArtBot. AI Horde, Stable Horde, Kobold Horde, Haidra are all the same thing (yes, there's too many names).
how to close comfy UI, do i just close the browser and cmd?
I think you just close the terminal.
Hey, so I just learned about this software called glaze, which artists are using to "protect" their images from lora and ti style training. Its an adversarial network that makes the minimum amount of visual change to an image, while throwing ai networks off of recognizing what style that image is. the idea is that when prompting, styles typed in will not correspond to the image output.
does anyone familiar know if this works or not?
where do i put samplers i downloaded?
there is a folder
is it the one called "extra_samplers"?
honestly i dont know
C:\ComfyUI_windows_portable\ComfyUI\comfy\extra_samplers
I find that there are enough options out of the box

Where could I ask for a prompt for a specific photo?
How do you upload a sketch and make it create a good image from it? I uploaded a pic but all it does is create a completely different brand new image
Is the teaching atlas of problem-solving ideas and steps for generating comic style leetcode problems using LLM and Stablediffusion?
hello. just a quick question. in dreambooth, is steps per image more important than total steps? like if i did a 10 image dataset with 100 steps per image for 1000 total steps and im happy with the aesthetic. but i want more flexibility so I add 10 more images. at 100 steps per image thats now 2000 total steps. is that necessary for a similar level of quality or can i halve the steps per image and do 1000 total steps and get a good result?
I'm trying to upscale using img2img and I'm getting this error:
NansException: A tensor with all NaNs was produced in Unet. This could be either because there's not enough precision to represent the picture, or because your video card does not support half type. Try setting the "Upcast cross attention layer to float32" option in Settings > Stable Diffusion or using the --no-half commandline argument to fix this. Use --disable-nan-check commandline argument to disable this check.
I already tried activating the "upcast cross attention layer to float32" option but I still get this error. I'm running an NVIDIA GeForce GTX 1660 Super. Anyone have a clue what might be causing this, and how to fix it?
Hey, what's inside your webui-user.bat ?
I'm not sure what you mean 😅 I use Stability Matrix. This is the full error log. https://i.imgur.com/5KjfHfU.png
and these are my img2img settings: https://i.imgur.com/r0HxqBE.png
I found the file you mentioned, webui-user.bat
It had the following:
@echo off
set PYTHON=
set GIT=
set VENV_DIR=
set COMMANDLINE_ARGS=call webui.bat
Okay at the line Commandline_ARGS=
you need to add:
--xformers --medvram --no-half --precision full
ok, do I disable the "Upcast cross attention layer to float32" option in settings?
Yes
ok I disabled it, saved settings, tried again, but I still got the same error 😦
hlo what is black theme name in automatic 1111
Can you post your img2img settings in #🤝|tech-support ?
yep, screenshot?
hi everybody, i have a question I am hoping somebody would have some ideas for.
i am using text to image and am feeding a summarized version of text through to be the prompt used to generate the image. The images being generated im using 50 sample steps and they seem to produce a nice job however the text being printed in the images are very difficult to read, even gibberish for most cases. Is there a way to make the text being generated within the images more clear? are we not prompting it someway that we should be so that text can be read clearly?
here is my example. hope somebody can help?
Unfortunately hands & words are the hardest challenges for any AI at the moment. I think that you should edit the words by yourself and add it into your image instead of waiting for SD to generate.
okay, well i wont be able to do that in my case as I am generating images on the fly and wont have time for manual intervention on the images.
sometimes it gets really close to producing excellent results and other times all the text is gibberish.
Does anyone use comfy ui, ? Need some help with installing the reactor node, it's frying Mt brain
Does anyone have a good doc for training self portraits with Lora or another way?
Is am1111 easy to install? I love comfy ui, but I'm wanting to use face swap
hello everyone,I'm working full-time on the open-source project SolidUI, which can generate any graph with just one sentence. I welcome friends to participate and contribute together.
You can also click on my avatar; there's a join link inside.
Project address: https://github.com/CloudOrc/SolidUI
welcome!!!
Does anyone know where the other buttons have disappeared? ||https://cdn.discordapp.com/attachments/1063492199821623296/1153365275824099368/Screenshot_2023-09-18_181947.png||
I'm struggling so hard to render a robot without a head, I tried "mech", "robot", "automaton", various "headless", and negative "head" and "eyes" in prompt...etc but it's rare I can ever get a robot without a head. You guys have any idea?
"mech" gives the most reliable results without a head but then it makes for a very bulky torso, which I haven't been able to prompt out
try vice versa big head small torso gives the mech more stability and movment
I'm not sure I follow, what do you mean?
I've also played around a lot tried to get a "slim" "agile" "runner build"...etc
like if the torso is small the rest of the body will be able to move more easily and a bigger head of ur chose will give the mech balance u know
unless u dont want a head
Yeah my struggle is that I don't want a head
then remove the head and keep the torso small big torsos make a mech to unstable the mech could actually crush its self if moved in the wrong derection
What do you mean "remove the head" 😆 I've tried all manners of negative "top heavy"...etc aswell
i meant remove the head from the picture i mention head earlier
i found something in my past, but not sure it is what you want + prompts seems to be not for headless robots.
I'm trying to generate a fantasy magically-animated robot made of wood actually, but yeah I understand the logic of what I want, I'm just struggling to hint SDXL to it
Sure, shoot 🤷♂️
i cant here, image is in general with images
Do you want large or small?
Have you tried headless?
i love talking and working on robotics
Yeah I've tried a lot of options, my best guess is that "mech" or "robot" aren't the best base
robot and pistol is only international words with our origin 🙂
true
I've played around with a few Discord and web image generators which say they use SD as their engine, and got fairly good results. So I installed SD itself on my Win11 desktop today.
So I had a look around the web for any articles to guide a total noob like me - but pasting prompts from those articles into my SD returned images that were nothing at all like the ones in those articles. I'm doing something wrong, or not understanding something.
Can anyone suggest anywhere I can look for patient guidance for a beginner on how to use SD properly?
you should have same model, same seed, and dont forgot to have vae. Can you post your pictures? @torn mist ?
what model are you using? 1.5, 2+ or SDXL?
and what UI are you using? A1111 or Comfui? there are other, but those are most common
https://stable-diffusion-art.com/beginners-guide/
may be helpful @torn mist
Thanks @plucky shoal . Sorry to display the depths of my ignorance, but how do I find out what model I'm using? Or which UI (I did see AUTOMATIC1111 flash by somewhere)
A1111 probably you are using. Model is in left top corner if A1111
@torn mist most probably default 1,5 version
Here's an example screenshot - https://imgur.com/KQeMjoP
Its a very simple promot, asking fo a mouse (I thought) but no mouse in the image. This actually one of the better pictures tbh.
Is that UI A1111?
you have not any model (checkpoint)
OK... how do I get one?
but it shouldnt work at all
(I feel so dumb, sorry)
post your pictures in general with images please. Hope it is o.k.
@torn mist here what is in folder stable diffusion
yes image you posted, folders, what is in folder Stable diffusion there.
Here cant be posted images.
Going to continue in "general with images" ...
yes
what is the purpose of saving grids? so you can look back and see how your prompts affected the images and such?
whats the newest stable diffusion?
#✨|sdxl 👏
Not sure if this belongs here or on another channel. #✍🏼|rules-and-tos says nothing about self promotion being not allowed on this server so I hope this is fine. If not, I'll move this message elsewhere.
I created a storybook generator that uses ChatGPT and stable diffusion. I created a survey to evaluate the system. You'll look at the images generated for 10 different stories and vote yes/no for if the images have a consistent style or not.
Shouldn't take longer than 10 minutes, will probably take less time. Here's the survey - https://forms.gle/ddPz2hje29G24Jz68
Hello guys, does anyone know how I can give the instruction to write a specific text in an image creation? I have already spent many credits trying without results
What do you mean? Do you want a sign with text? A paper with text? A shirt with text?
Thanks for responding, yes it is a banner style logo. I have tried something like: description, background, logo, with text "x". and complies with everything except the text
What scheduler are you using?
what do you mean? sorry, i am new on this
Free
What chat to share models?
Guys, whats is ur fav/best anime model for 1.5/SDXL?
Holoanime for sdxl is quite nice
I personally quite like Counterfeit for 1.5
yall ever heard of tuneflow the melody AI generator?
Hey Guys,
I'm wanting to train a model on a specific pixar-like character. I've generated 20 images of the sam(ish) character and I want to be able to prompt that character via dreambooth or something similar. Any tips for how I'd be able to do that?
what's the announcement
its not an announcement. It's just pitching stable audio again
exact same annoucnement as last week, just fresh on the notifications . if that's how @karmic brook going to do the annoucnements page i'll mute em
it takes all my effort to fight notification spam
Considering the amount of people using depthmaps with ControlNet Depth, is there any way to do the opposite, and generate a depth map from a render?
or to generate depthmaps in the first place?
any 3d renderer can output a depth map of the scene
by "render" I mean generation, not a 3D render. I don't have the 3D information, it's what I'm looking to extract 🙂
any easy installations for local use for StableAudio? I have only 6GB vram tho but maybe vram isn't that important since we are not talking about images
you can generate a depth map out of a 2D image or a 3D image doesn't matter
My question is how lmao
But nuuideas was saying that if I had a 3D model I could just render a perfect depthmap directly, which is fair, me saying "render" wasn't exactly clear
the depth extension from auto1111 is one way to do it
a1111 doesn't work for me for some reason, is there something similar for ComfyUI?
I was hoping for something based on image diffusion, I know there's a lot a tools out of there that can do a decent job (even Davinci Resolve has a node), but considering the amount of depth -> img generations I was hoping something trained a model for that
there was this opensource program that also did it although idk if it was opensource it was free anyways where you also can generate displacements maps etc etc
Materialize?
oh ye could be
Tool looks nice but is mostly for texture. I've used tools like that in the past but they don't really behave so well
ye, but a depth map is basically a height map same data
I know, I'm not looking to use that for flat textures though
a height map creates depth
you can basically create a 3D mesh with them like those plane things that you use for 3D printers
Right, depthmaps tend to not (always) be orthogonal projections though, and they're rarely linear, but I'm saying something like Materialize couldn't help me because it's not made to generate depthmaps from random images, it's made to generate heightmaps for diffuse textures, normals...etc
And as much as I like elegant algorithmic solutions, I don't think the tool is nearly as good as what image diffusion could do
Hello everyone, new SD noob here, I need some assistance regarding controlNet, I keep getting this same msg every time I try to generate a new image, "RuntimeError: mat1 and mat2 shapes cannot be multiplied". If anyone got any spare time to explain how to resolve this issue
Has anyone tinkered with extensions through stable diffusion api?
I'm trying to use root extension through api with "alwayson_scripts" field but there is no documentation around it
anyone know why my automatic1111 reloads the model every time I change my prompt?
extremely frustrating
@sharp panther it is because you are using things for some version for different version.
SD is so dead right now
I keep getting the not able to finish the dream error
as a doorknell
Hello! Anyone willing to help an SD newbie out?
Have only been dabbling with midjourney and now venturing into SD to make my generations "pop" even more. I'm trying to upscale and sharpen my image via "img2img" with the "Ultimate SD upscale" script.
I've tried different settings but however I go about it I just end up with an image that is less sharp than the original
im trying to use controlnet with a 3080ti but it takes 20 minutes per render, am I doing something wrong or is that how long it takes for everyone? using A1111
It should be much faster
What are your settings?
huch? Why?
hello
I am new at stable diffusion I was using stable diffusion 3 month ago - I had only 4 gb of vram back them but I have upgraded my pc
I have purchased rtx 4090 to learn all the concepts of stable diffusion
I have installed auto 1111
and I have also installed sdxl model
I don't know how to setup refiner and I have also added mikubill controlnet extension and downloaded depth - canny - openpose models for controlnet
how to activate it?
you have to download the refiner model, it is a separate model
then it is loaded as a regular model, and you can generate. But it is best use (I guess) if you do it in the generation process, right bellow the Sampling steps.
so it switches between the base model and the refiner, you can the % of steps each one is going to generate
I have downloaded refiner model
now the problem is how can I have installed extension of controlnet but its not showing
😦
Hey, I'm using Automatic1111 now, other webuis I've used before had previews of styles (like "photorealism" or some artist) and when clicking the preview the style would automatically be added to my prompt. Does something like this also exist for automatic111?
Hi everyone,
I'm looking for someone who can help me with AI art. I've tried to get into it myself, but I work a full-time job and film and edit videos for YouTube and I am a single dad, so I just don't have the time to dedicate to learning everything I need to know. I also have some hardware limitations.
I'm looking for someone who can generate realistic art and who is familiar with ControllNet. If you're interested in helping me out, please send me a message.
In return, I can offer to promote your social links on my YouTube channel and in my Discord server. I have a fairly active community, so this would be a great opportunity for you to get your work seen by more people.
Thanks!
no sdxl controlnet
do you know if SDXL is trained with names of plants, species and such? I tried to generate "empress tree", which is a name for a tree with purple flowers. It instead generated a woman with empress clothes in all green fashion
Oh
I mean people would use SD without controlnet as well.
How do i generate pictures?
Where can I find the best tutorial? For how many images I need to create a character in SDXL? What's an ideal number of images to use for training?
Hi, is stable diffusion banned on google colab? It keeps disconnecting me while using it, is there like a way to stop it from doing that?
I click the image preview while its being drawn but it switches to the previous image.
Is there a way to get it to preview the current image even though its not complete?
wow dalle3 is amazing. is it game over for sdxl? will stable diffusion ever recover from this? https://x.com/willdepue/status/1704558773216661572?s=46&t=t7Hr_X6l1AqBhbnh662qjQ
it does impeccable text too
it looks promising, but lets see 😄
the old internet of yore coming back to haunt us
it was supposed to be a paypal competitor and they got acquired. it's how elon got bootstrapped. Now he's rebranded twitter as x and resurrected it
dalle 3 looks like an iterative improvement but nothing "amazing". people hyping it's pixel art capabilities are blind. this was called shockingly good pixel art... https://twitter.com/willdepue/status/1704560786121834768/photo/1
most of it is just noise
I convinced that some new update around 1.6 made webui a large amount slower... It went from 5min per gen to around 8
likely a localized issue or you're skewing your own memory
5min per gen tells me theres issues right out of the gate. that sounds like CPU mode
Using a 1660ti and outputting 512x768 with 2x hires fix should not take 8 minutes
but lets see it in practice
Is dpm 2s a karras slower than other ones?
if that's real world usaged. Remember when videogames surged from a small to ALL ENCOMPASSING market? we got a decade worth of cinematics for commercials instead of actual gameplay
I moreso meant slower than the other dpm++ samplers
openAI are likely staging these conversations. Training specifically for them. Hand crafting larry before feeding it into the training data. Trying 100 different reconstructions before they took the clips that looked like a natural conversation. Who knows how they actually did it. It's such a highly edited video that it's certain it's not real UI use
yeah, also it wont be easy for them against Midjourney, Stable Diffusion and Firefly
but it can eventually accompany one of those
since people already pay for ChatGPT Plus
openAI are partnered with Microsoft. This is just a marketing blitz. They aren't worried about MJ or SD
or use it otherwise
All windows 11 have chatgpt licensed now. They got their funding secured
well i guess they dont sh*t in their pants because of the competition in AI art area
chatgpt api access is more about teaching the market these tools can be used. It's a halo product.
you mean the Bing thing?
it's their biggest success
actually i forgot that Adobe and Google partnered as well and Bard will have the same thing as OpenAI is doing with DALL E 3 now but in their case it will be with Firefly
and Bard is free...for now
bard isn't in canada since google hates canada's new laws
yeah they're pissed about our new media laws so they're holding back services
interesting
hey, is there a good website that showcases some great stable diffusion made images with their models and prompts?
$100 a month makes it a huge scam. It doesn't even implement inpainting very well. I can tell there's no controlnet models at all.
and that camera angle control? 100% fake controls. Your service probably signs up for auto renewel and is hard to cancel too
@copper crystali mean even if not a scam, the price tag is bigger than Adobe CC all apps plan
this one is supposed to be alternative to that
maybe it draws credits from GPT-3.5 or 4 or Dreamstudio API
Yes, but we had it in July. We've not had it for three months, and there's no sign of a return.
SD went backwards
I am so unsure if i should just say F it and get a budget nvme, seeing as even comfyui doesn't load models faster than 1.3GB's on a gen 3 that can read at 3GB's anyways, or get a "more durable gen 4".
why did it disappear?
I don't know. When SDXL was released, there was no controlnet for it. Now there still isn't.
NEver mind. Didn't read the damn SDXL part 
there is a good reason why theres no controlnet for sdxl but luckily we have controlnet loras, thats all we will get for now
for some functions you dont need controlnet anyway tho if i remember correctly
although im not sure since i used SDXL for 2 weeks before i deinstalled
there are
diffusers team released a bunch
just the same canny depth
how you mean the same? you said none a second ago
more of the same as controlloras depth,canny,inpaiting
but that isn't none?
"luckily we have controlnet loras, thats all we will get for now" << why say wrong things then argue that they're right things?
oh yea u right luckily we have an english bot here to correct when something goes wrong
hmm
people sure do love to spread misinformation and lies. i wonder what that's all about
not everyone can be as smart and informed on current topics as you,as humans we sometimes miss updates on current themes
but, you knew about the diffusers controlnets. you told me what they were..... holy moly what even is this
sounds like a redditor thread would you like to continue i can switch to a more agressive LLM if you want
before i came in everyone was talking about how controlnet is non existent. before i could even finish mentioning which ones diffusers had released, you told me that doesn't count
there was a team effort saying that controlnets don't exist here.
hi! question for any software developers, m2 macbook air 15 or m2 macbook pro 14... and is 16gb ram enough to not feel strained
This isn't the first time i've come in and said "y'all are wrong and this is why" and people get mad at me and act like victims because english isn't the primary language.
passive aggression is aggression and always rooted in anger
i wouldnt be the person i am without you
being wrong is a good thing. it's your opportunity to be right next time. Don't be so sulky when you find out how wrong you are. instead celebrate
of course we celebrate,we have you
so this means we're still stuck on the same controlora level of releases cany depth and inpaiting?
aye ma'am i see a recent openpose release too u tried that?
i get errors with openpose sdxl ill try to ask in support chat 😊
You don't need it. But for getting a pose or consistent style/look, it is much better than native model.
Without it, everything must be made into a LORA
There is no cany for SDXL, at least not in SDNEXT
literally just covered that there is controlnet models for sdxl
did some youtuber declare that it will never happen or something? People keep coming up with this idea that it's non existent for sdxl. i see it over and over
thats y i asked here if there was more other than the same old 3 we already know but ty 👍
https://huggingface.co/lllyasviel/sd_control_collection/tree/main all of these work
all the ones marked sdxl
Those work for SDXL in which UI's?
Hello, i'm installing stable diffusion on local and i have some questions, which channel should i ask?
anywhere that controlnet is available more or less. comfyui or the extension for auto both support all those. check the model card. all the info is there
update sdnext. it just uses the contorlnet extension for it's purpose
@frail prism hi, thank you for your stickers workflow! I have one question, as iam not good at workflow. Whats are parameters to get out of aliased outlines? I played with it a bit but cant get good result. Played with feather and so.
I'm not sure what means "aliased results". Stickers without borders?
borders bud jagged? if proper word in english
mmnt
Try to shortcut "Mask Smooth Region" node. Border will become more jagged.
i want not jagged, I mean antialiased i want 🙂
@frail prism i think it is because rmbg cropping it not smoothly
Oh, well, increase sigma parameter in mentioned "Mask Smooth Region" node. The mask will become smoother.
@plucky shoal Or you can draw mask manually in Mask Editor and use lower nodes line (with blue node at the end).
thank you! @frail prism will try!
@plucky shoal I can't paste image here, but note that there are two "Mask Smooth Region" nodes. One used for automatic masks (higher), and one for manual masks (lower).
And to use manual mask, connect blue node "Mask Erode Region" to blue node "Convert Mask to Image".
o.k. i played only with bottom i think
Hey there
Hope you are doing well
I am a senior full stack developer who has full experience in web and AI development
So if you have some projects, please let me know
thank you
https://github.com/typhon0130
https://figma.com/@typhon0130
Hey, you can ask in #🤝|tech-support for installation help
guys what are the dimensions for sd1.5 models
SD 1.5 models got trained on 512x512 resolution
dear fellas, u guys know a free method to geneare music video?
can i get 16;9 images using 1.5?
can someone explain the spaghetti thing to me?
Yes you can by using a resolution like 540x960 and then upscale by 2
okay thanks
I had to apply --no-half to my startup to get SD to do sketch inpaint. It's inredibly slow. Is it generally that much slower than the other operations or is it related to the no-half thing?
I wonder why when I generate image about hald body portrait...
it look beautiful
but when I try to generate full body
It's suck, ugly as hell
The details are a lower resolution. This affects two things: one is that the lower resolution already looks bad. The other is that the images are generated in little chunks that are lower in resolution than the final pixel resolution and they guess what the pixels in between are supposed to look like. This gets very hairy at low resolutions.
Good afternoon, everyone!
hello everyone,I'm working full-time on the open-source project SolidUI, which can generate any graph with just one sentence. I welcome friends to participate and contribute together.
You can also click on my avatar; there's a join link inside.
welcome!!

Upscale your videos to 4k on free google colab: https://github.com/yuvraj108c/4k-video-upscaler-colab
with real esrgan... 🙄
would be better off with a nearest neighbor algorithm for video
Anyone having issues with using Stableaudio? It's not letting me login, the "try it out" button does nothing. :/
i think this is common as service is very used. Try it and it let you login
are there any prompts to generate high quality vectors?
There a furry model that does good fur? Most furry models I find lean more towards humanoid instead of fur
I like the look of soft plush fur
Fur you can run your hands through
Like they use conditioner
Like very fluffy
Nice fur
Fur
I love you
Be my wives husband please
Ive tried yiffymix and that doesnt do good fur, it is just humanoid with a "x" animal head
Kemonofurry model does better a bit, it does nice tails
The model I prefer is EasyFluff but it's a vpred model. It would require cfg rescale.
hello guys, do I need sdxl vae when dreambooth training?
for sdxl
as an option to kohya_ss, I believe I dont need but I ask here just to be sure
Whats a vpred model?
https://www.reddit.com/r/StableDiffusion/comments/13joe98/sds_noise_schedule_is_flawed_this_new_paper/
https://arxiv.org/abs/2305.08891
It's a model that has a different loss training schedule.
Interestinv
Made a bat file for running comfy in one click from desktop for those who wants it
Just paste in notepad, save as "comfyui.bat" for example and run.
replace "user" with what your pc's user name is.
Replace "G:\Stablediff\Comfyuimanual\ComfyUI" with the path to your comfyui's root folder
set conda_path=%windir%\System32\cmd.exe "/K" C:\Users\user\AppData\Local\miniconda3\Scripts\activate.bat C:\Users\user\AppData\Local\miniconda3
rem Change directory to "G:\Stablediff\Comfyuimanual\ComfyUI"
cd /d G:\Stablediff\Comfyuimanual\ComfyUI
rem Run "main.py"
python main.py
do you think controlnet will ever release scribble for SDXL?
there's a sketch control-lora for sdxl @ https://huggingface.co/stabilityai/control-lora which is pretty close in concept?
hi anybody has ever trained stable diffusion for video generation?
Any idea when we will have any big improvements in prompt following in open source models? Dalle-3 is insanely good at it
Hello everyone. I'm trying to create a set of 2D monsters for an indie RPG video game project. Non-animated, 2D full frontal view of monsters (oldschool 🙂 ). I just started using SD and noticed that what is probably currently preventing me from achieving that with any kind of consistency is the model i've tried using (sd 1.5 and RV). Do you have any model in mind that is better at this task ? Any tips or things I should look into 🙂 ? Thank you in advance !
⌨️👁️
Hi all im using SDXL offline, i see there are bots here on discord are these bots generating images and are we aloud to use them?
hey, yes you can use the bots in the bot channels for free 🙂
hey sorry to bother but can we move here https://discord.com/channels/1002292111942635562/1002602742667280404 really quick, or do u not have time?
aku
sayang 👹 kamu
juga
sayang 
There's no proper video generation yet afaik. There's animatediff, but it can only do max 24 frames at 8 fps for a decent framerate, and even at 512x512, it will use ram as 24GB memory is eaten up.
there's workloads that is sorta tedious that can do longer, but needs batches of X amount of frames, and have to manyally be post processed per batch
I get this "RuntimeError: The size of tensor a (384) must match the size of tensor b (320) at non-singleton dimension 1" When trying to run the model "control-lora-sketch-rank256" with preprocesso "scribble_pidinet"
Do you know how to fix?
same with any preprocessor
Hmm, can't decide whether to try cheaper storage being Crucial P3 for 70 euro, or just say F it and get 970 evo plus on sale. for 94 euros. As it's only purpose will be to load models, so dram is not needed at all.
Did u find a way? It keeps disconnecting me too
hi
thanks..wondering how synthesia.io does it hmm
No there's just no way anymore, I gave up and started using it on my pc instead, other websites are way too restrictive to use for free
i found a very good picture from an AI Generated https://www.instagram.com/p/CxaeCGVNTk-/?img_index=1 but dont know how i can generate it and whit which promts or checkpoints...
is ther some way to tell if an AI image generator is diffusion based on the output?
Same, just started using it locally. And its much better
Ngl I liked google colab a lot more, it was way faster and it didn't have as many errors, but compared to the free ai art generators then hell yeah, it's not even close xD
How do I tell Discord not to notify me about stuff in the #1072240143521554592 ? I already set the notifications on it to "Nothing" but it still shows me a badge whenever they tag @challengers in there...
⌨️👁️
i need help...
venv "C:\Users\new user\Desktop\A1111\stable-diffusion-webui\venv\Scripts\Python.exe"
Python 3.10.6 (tags/v3.10.6:9c7b4bd, Aug 1 2022, 21:53:49) [MSC v.1932 64 bit (AMD64)]
Version: v1.6.0
Commit hash: 5ef669de080814067961f28357256e8fe27544f4
Launching Web UI with arguments: --xformers --autolaunch --medvram
Loading weights [5aacd1c8c1] from C:\Users\new user\Desktop\A1111\stable-diffusion-webui\models\Stable-diffusion\ParchartXL-1.2.1.safetensors
Running on local URL: http://127.0.0.1:7860
Creating model from config: C:\Users\new user\Desktop\A1111\stable-diffusion-webui\configs\v1-inference.yaml
To create a public link, set share=True in launch().
Startup time: 33.2s (prepare environment: 10.4s, import torch: 7.5s, import gradio: 3.5s, setup paths: 3.6s, initialize shared: 0.5s, other imports: 3.3s, setup codeformer: 0.2s, load scripts: 1.7s, create ui: 1.4s, gradio launch: 1.3s).
Applying attention optimization: xformers... done.
Model loaded in 16.6s (load weights from disk: 2.4s, create model: 2.5s, apply weights to model: 0.1s, apply half(): 3.8s, calculate empty prompt: 7.8s).
0%| | 0/16 [00:06<?, ?it/s]
*** Error completing request
*** Arguments: ('task(iq9fv8j43ns4ivq)', 0, '', '', [], <PIL.Image.Image image mode=RGBA size=820x616 at 0x2649ED15390>, None, None, None, None, None, None, 20, 'DPM++ 2M Karras', 4, 0, 1, 1, 1, 7, 1.5, 0.75, 0, 512, 512, 1, 0, 0, 32, 0, '', '', '', [], False, [], '', <gradio.routes.Request object at 0x000002649E8522F0>, 0, False, '', 0.8, -1, False, -1, 0, 0, 0, '* CFG Scale should be 2 or lower.', True, True, '', '', True, 50, True, 1, 0, False, 4, 0.5, 'Linear', 'None', '<p style="margin-bottom:0.75em">Recommended settings: Sampling Steps: 80-100, Sampler: Euler a, Denoising strength: 0.8</p>', 128, 8, ['left', 'right', 'up', 'down'], 1, 0.05, 128, 4, 0, ['left', 'right', 'up', 'down'], False, False, 'positive', 'comma', 0, F
i can't get stable diffsion to do anything. i have no idea what the problem it
come to #🤝|tech-support with info about your gpu and what your tried to get this error
how do I avoid generating images with the camera too close to the ground and out of focus background?
https://i0.wp.com/digital-photography-school.com/wp-content/uploads/2016/12/Camera-on-ground-perspective-02.jpg I want to avoid this angle which is often what I get
that type of error (tensor size mismatch) generally indicates a mixup of model files (eg non-SDXL model with the SDXL controlnet or something)
hey guys, is this server bot better than the stable difussion app itself?7
amazing prompts updated there is no such other resources : https://github.com/FurkanGozukara/Stable-Diffusion/blob/main/Amazing-Prompts-List-For-Stable-Diffusion.md
whats the differance between the different bots?
Very little, but bot 1 has controlnet available.
Hey I'm Dhruv and I would try my best to nurture the friendship with any of you.
My energies are always there for somebody's help, usually in python.
Hey chaps - wanna make basic 90s style rave vj loop style images. Like early Bryce 3d. Are there any models I could use!? Key prompts u could suggest?
hey im new to stable diffusion but ive been using MJ for a while. are there any outstanding fundamental differences in image generation between the two?
Like MJ is aesthetically pleasing and gets so many things right, all without having any understanding of what its creating or how anything is supposed to function. Does Stable Diffusion have the same issues?
In Stable Diffusion you have more control over the output as in MJ.
Especially when you install it localy on your machine you can do much more than in MJ.
In Stable Diffusion you can use different models to get different styles. You won't get any aesthetic perfect image from a 3 word prompt like in MJ.
But the SDXL model for example is pretty good with short prompts
guys can anyone tell me prompt for Guts from berserk lying in thunderstorm image expressing extreme sadness and solitude in manga styled artwork
Thats why some of them use MJ in combination with some other tools
But still you cant train your own models etc. Not yet at least
And now DALL-E 3 seems to become a serious threat for MJ, for SD not soo much because many in the SD community prefer some of the SD advantages
Hi guys, I updated automatic 1111 on my mac, and now it won't recognise my .safetensors models. can anyone please point me in the right direction to solve this. I've tried downgrading python
you can try to delete the venv folder and rerun the webui
Anyone installed the new nvidia driver?
Does it still negatively impact the performance of SD?
I had to roll back to the one in May
hi
Hi guys, I have been using some LORA models for certain specific actions coupled with LORA models for celeberty faces. However, I am kinda getting the worst of both worlds, the ressemblance is a lot less and the pose always turned out to be kinda just standing still and looking at camera. I have tried messing around with different weight and the result is similar. Is there any way to make the 2 LORA work better together?
So I’m getting into machine learning in my new job at work which led me to finding stable diffusion. Setup on my amd 6700 XT gaming rig last night but as with gaming, I see the issue of not having tensor cores , I’m a hardcore amd junkie but if I am getting into ML and SD , I assume it’s time to ditch the 6700 XT ? If I built a cheap rig just for ml sd what’s the cheapest nvidia card to go with? And will this work with dual gpu’s?
Think you, will try this
damn, no luck. it recognises the base model that comes with Auto, but no other models in the fodler. 😦
Tell me prompt to generate strawberry elephant anoano wilkum whatsapp status
Does it not recognise alle safetensors or just some of them ?
discarding roop, with what other tool can I do face swap on A1111?, to be able to do something like remini
what's the best-best gpu for stable diffusion?
1min
Im completely new to using ai image generation. I have low resolution images that i want to upscale. Ive heard about img2img, controlnet, and real-esrgan but i dont really understand what they are and they relevance they play into upscaling. Also, does the input image resolution affect the quality of the final product (i have the same image as 640x480 and also as 960x720). In my head, it would make sense that the latter would be better to upscale with but idk if these programs work better with smaller images or not
What is your most used Checkpoint? Realistic Vision, epiCRealism, epiCPhotoGasm, epiCDream for best creations?
Dreamshaper is still pretty good.
🤔 I'm gonna use random messages here as prompts
How do people go about making those loras for the world morphs?
Welp. Solid pdf prompt book....for the rookies 😉
I need to know how I can engage two different characters in one image doing romance?
How to make stable diffusion work again after copy pasting the folder on another pc?
Is there any extension i can use in automatic1111 where I can engage two characters in one image with two different seeds and prompts?
All unfoirtunatley. There are about 8 or more in the folder. but only recognises the one that came with the install
does not recognise all I mean, except the one that came with the install, which is weird.
How big are the safetensor files ?
from 2gb to 7gb
can you make screen with directory where models are? @worthy lichen curious, what can be problem. Only i can imagine suffix is somehow wrong...
Yeah sure. I was trying to put in a screen shot but don't think I can put images here.
Directory is: Automatic1111/stable-diffusion-webui/models/Stable-diffusion
put it in that chatt bellow
The weird thing is that it was working, and then after update just didn't recognise them.
seems o.k. i just have / other way round
probably some extension, have you not some sort of civitai loras or similar?
yeha thats a mac thing
pictures in general with images
no. Didn;t download Loras yet. I did get another ckpt just to see if it was a safetensors prob, but doesn't see that either. So frustrating 😦
Hello guys, I've searched the conversation but I have not found a solution to the issue I encountered after installing SDXL, It works to some degree but my controlnet seems broken 9/10 times.
this is what I have in the log:
*** Error running postprocess: D:\AIstuff\Image\Automatic1111\New folder\stable-diffusion-webui\extensions\sd-webui-refiner\scripts\refiner.py███████████████| 20/20 [00:12<00:00, 1.48it/s]
Traceback (most recent call last):
File "D:\AIstuff\Image\Automatic1111\New folder\stable-diffusion-webui\modules\scripts.py", line 651, in postprocess
script.postprocess(p, processed, *script_args)
TypeError: Refiner.postprocess() missing 3 required positional arguments: 'enable', 'checkpoint', and 'steps'
yeah. Reinstall a couple times. Thanks so much for your help. I rented an online Auto install. Graviti. Seems to work well. And using invoke locally. bought this mac before the AI boom, and now really wish i had a PC
@worthy lichen
can you check this if related?
https://github.com/huggingface/safetensors/issues/330
if related there seems to be something to try
Awesome, thank you so much. Going to read through carefully. Does seem related.
👍
I am using automatic1111
would try to disable refiner @quiet kernel but i got realy bad ideas
thanks I'll try righ now
@worthy lichen it seems updating python to 3.11 which on windows wont work is one of solution
Chat bots are literally better than Discord
won't work on windows?
same error after unchecking refiner
on windows only python 3.10.6 -> 3.10.11 or so
I wonder if I should just reinstall from scratch at this point
You can, shouldnt be hard. Or try delete only VENV folder and run webui user.bat again
ok, trying
Edit the webui-user.bat and at the line Commandline_ARGS= you add: --device-id 1
Cool thanks. Yeah I had 3.11.something, downgraded back to 3.10.6. but when AUTO Launches, says its launching 3.10.13. was digging late til 2 am on friday trying to find out how to fix.
and other things do not helps as well?
Doesn't seem like it. deleted the venv. So weird. Is there an extension for importing models using the gui. maybe i google that
been looking through all the resource stuff for hours, and i guess im just dumb or blind, but i cannot find any documentation or instructions on how to download and install stable diffusion on my own PC, only source code from the github-alike and how to run it from a cloud service.
can someone please point me to a link or a set of instructions on how to just download and install this onto my own PC?
really wanna try this
there are quite a few guides. It is a bit manual.
I'm testing this since I have probs with my mac installation: https://www.diffus.graviti.com/
its paid tho
@worthy lichen i think it is easy to follow. It is probably best how to install it.
Thanks, I agree. PS found a backup of old Auto install. and it works. Damn. wasn't as stupid as i thought.
will stable audio be installable on your own pc like SD?
Now I broke Invoke AI. yay
hi guys
i need help finding R-ESRGAN-4x+ upscaler, where can i get it? i searched everywhere
oh man thank you, this is exactly what i needed! i think!
Does comfyui give better results than a1111 whaddayathink
I use stable diffusion A1111 …. Are there any other free app that are as good as SD? I think Dalle and clipdrop are both paid, right?
tf
I did just that it is reinstalling righ now
this is where I'm at now:
ule.py", line 1501, in _call_impl
return forward_call(*args, **kwargs)
File "D:\AIstuff\Image\Automatic1111\New folder\stable-diffusion-webui\modules\sd_hijack_optimizations.py", line 486, in xformers_attention_forward
k_in = self.to_k(context_k)
File "D:\AIstuff\Image\Automatic1111\New folder\stable-diffusion-webui\venv\lib\site-packages\torch\nn\modules\module.py", line 1501, in _call_impl
return forward_call(*args, **kwargs)
File "D:\AIstuff\Image\Automatic1111\New folder\stable-diffusion-webui\extensions-builtin\Lora\networks.py", line 429, in network_Linear_forward
return originals.Linear_forward(self, input)
File "D:\AIstuff\Image\Automatic1111\New folder\stable-diffusion-webui\venv\lib\site-packages\torch\nn\modules\linear.py", line 114, in forward
return F.linear(input, self.weight, self.bias)
RuntimeError: mat1 and mat2 shapes cannot be multiplied (154x1280 and 768x320)
can't copy the whole log
Can someone eli5 about the hi-res fix? I find text about it everywhere but not what files i should download or how to get it to show up as the upscaler version some people have in their A1111's UI....
Hey, friends,
I've been playing around with training Lora models for a long time, so I decided to get some experience training Dreambooth models. As a test, I created a model from about 500 images, set to 42 steps and 10 epochs, with a learning rate of 0.000005, with the intention of teaching it more images and styles later... Unfortunately, I found that in Kohya_ss for Dreambooth training, the option to build on previous training which is present in Lora training is completely missing.
So my question is, how can I train the model further? Should I select it as a source model and train new images on it? (if so, should I include the original images or just the new images?) or do I just have to train a completely new model, leave sdxl_base as the source model and then merge these models together or is there another way?
Thank you in advance for the answer!
is there a way to get a heatmap of what parts of the image a lora is affecting the most?
What are the best free websites for cloud models
Usually I just switch between Prodia and Patience.AI
Prodia for quick gens and Patience when I want to use SDXL or more detailed art
I should add I (still) can't use a local model because I'm on Chromebook unless there is some easy workaround for this
there is no extra "continue training a checkpoint"-option because choosing your source model is always continuing training it 😉
There are a lot of potential causes for that, but to make things easier on yourself you may just want to use one of the installers from https://github.com/brkirch/stable-diffusion-webui/releases to create a new install. If you are using an Apple Silicon Mac then you can try out the installer for the latest experimental version (you can download the application directly too if you'd like but currently I'd recommend the installer as it directs you to install your SD models first - the application currently downloads the SDXL model without prompting on launch if no models exist in ~/Documents/Stable Diffusion Web UI/models/Stable-diffusion). If you're using an Intel Mac then try the v1.6.0 release (https://github.com/brkirch/stable-diffusion-webui/releases/tag/v1.6.0-mac-builds) for x64. Either way if you have any problems loading the safetensors models with the install generated by the installer then please create an issue and I'll take a look.
since it already is trained on your first run of images it only needs the new ones. Don't train a new sdxl base and merge afterwards. Because then it can't really built up on each other
So I’m going to spin up an aws g4 instance to see if a T4 will work better than my amd 6700 XT . I put Ubuntu on it. So what’s easiest way to get SD up and running on it?
Else I’m considering replacing the 6700xt with a 3060 12gb . Getting as working on amd has been a pain
Do you guys know how you call a LoRa model in Fooocus? Is it the same as in A1111 with the trigger word and the model name in <______:1.0>?
Thanks a lot. Will do this
Thank you very much, I was worried not to overwrite the data learned in the first training in this way. 🙂
this is what i was telling people about heygen they werent beliving me : https://twitter.com/GozukaraFurkan/status/1706031986941268327
What's your issue with the 6700xt? you may need just the right cmd args in your webui-user.bat
Python not found fix?
Install python with "add python to path" checked.
If you already installed python then you can run the installer again. Click modify and then select "add python to system variables"
After that delete the venv folder and rerun the webui-user.bat
i have a question, if i want to generate an image that contains two trained subject, how would i do this? Lets say i have a model for each trained subject, can i merge them both? How would i merge these 2 models?
You can merge them in automatic1111 with the checkpoint merger tab
Hey
I have a question:
I wanted to try the new Toonyou Gif Animation
But everytime i try to generate i get this error
RuntimeError: CUDA error: invalid configuration argument Compile with TORCH_USE_CUDA_DSA to enable device-side assertions.
Can anyone help me with that ?
hello
has anyone successfully installed stable diffusion on macbook pro?
i am getting some errors please help
Is it better to make a lora or a lyco 😮
You mean AnimatedDiff ?
Then yea thats a common bug. Its fixable by go into the AnimatedDiff settings and change performance to Torch 2.0 instead of xformers.
But then it will be slow.
Hopefully an update will fix this
Thank you it worked
But I realised I don't have enough VRAM for this process
do i need to use prompts if i have a model trained on my own paintings, and i want to give my 3d renders a paint filter with image2image? ive seen some tutorials and i see people use prompts both when generating with image2image and when they train the model by describing every image
so if the 3d renders are very similar to the result i want, do i need the descriptive prompts of elements in the image except the style
Hey Gs, Which files should I download along with A1111?
@warm junco
what is the recommended to get higher resolution and video to video, can you help me plz
For the installation or usage?
The way for highres is upscaling
Is it very hard to make a good lora/lyco...? 😮 I wanna make 1 of a celebrity that isn't covered but I'm anxious, it looks tough xD
Please I can’t find the site link please send me
anyone have any idea why im getting like bright green accents in my img2img when i didn't prompt that color and that color is not in the original image. every model is doing it
can you post example? In channel below?
the green doesn't show up in the preview. the previews look nice and then when they're fully done the green appears
o.k. seems VAE issue if on end
can you turn vaes on and off
what model are you using 1.5 2.1 or SDXL?
1.5
what UI are you using? Vae is important
a1111
go into settings, there show all, and search for quicksettings. Here choose Vae, it will display it in top part of UI and you will be able to switch them or use none or Automatic, which is probably best
is there a vae in there by default, or did i dl something?
it is called sd vae there
it was on automatic though anyway
it is usualy baked in model. You can choose with baked vae and use automatic.
or you can DL if not wrong, size should be about 300MB?
download what
vae-ft-mse-840000-ema-pruned.safetensors
Try find this file and put it in vae folder, if not wrong
ye thats the only one i have
I'm also interested in VAE's. There's an "Automatic" setting - is this generally a good idea or no?
im using pure white backgrounds. idk if that could have anything to do with it
Hi all, any AI artist or just artist in general here? I would like to gather some suggestion on an AI art platform mobile application that im building. I cant really go into much detail about the business model but its basically a platform where u create art for free, follow ur favorite artist and like their artwork. Free to create art and video and some other thing.
it is good idea.
If it not happen in preview, and in last step, it is very probably VAE
Please DM me, i really want to build a platform for any “artist” to freely unleash their creativity via art and connect and support with each other so i would really love to gather what artists truly need
should i have the base model 1.5 checkpoint even if i only use other models, or does it not matter?
i dont understand now. Vae should be per model
1.5 for 1.5 SDXL for SDXL and so on.
if you asked about vae model, i would switch it to auto
idk if you were responding to me, but i wasn't talking about vaes. was kind of a side question
If I understand correctly, you can only use one checkpoint at a time unless you are also doing refining - then you can use two - basically
(A1111)
Question about CFG scale:
How come some renders become absolute trash if the scale is higher than "what it wants"??
Do we have something like Shuffle for SDXL ControlNet?
The reason I ask, I'm trying to figure out a way to generate an image using a very specific color. Do you know of anyway I could make that work? I have a image of just the color that I was trying with different ControlNets, but to no avail. Tried Reference, Recolor, and IP-Adapter
how to install stable difussion xl on linux? pls ping
is there a tool that automatically inpaints a whole area like a whole white background?
that's interesting idea, huh
i installed inpaintanything extension but cannot get it to show up as a tab on my ui
yo
Tried alot get about 2-3 it/s, I spun up a aws server with Tesla t4 and getting twice that, cuda for the win
How hard is sd on hard drive bandwidth? Does having a fast ssd make a difference? I assume it doesn’t cause no way sd can generate data thT fast
loading models is a lot faster on a good ssd
since they are pretty big files
The aws server has a 125gb nvme, but it’s not persistent, aka if you shut down you loose the data, so either I have to sync it to slower storage and script a restore it ot just get some standard ssd added on
ngl i have no idea what youre talking about sorry
can you use a pen tablet with A1111 ? to make masks in the inpainting tab for example
That sounds cool
if it works, that would be awsome
I paint anything add on does auto masks not tried it yet
But don’t listen to me I’m a noob ha ha
how do i get stability ai staff role and can i do infinite generations with it
who do i gotta ping 😵💫
Would like to play around with stable diffusion, controlnet and qr monster to see what i can mimic from reddit. Is there any... guide ? I guess ? to follow so that I dont install it incorrectly? I guess Im asking like, I am not sure if i need to mess with automatic1111 for that, or not... or what method I need to go about installing it, if that is my end result need
yes
that is awsome. thanks
I was doing it today
can anyone help? im getting these out of place bright green highlights on my img2img gens, where there is no green or anything in my prompt or in my original image.
it also looks like the highlights turned yellow/gold in another model
it will be like a thin spray-like mark around the edges of some objects
Is there a Discord app that doesn't suck
Official one is super laggy on my device
can you share an example? also, are you using more than 1 lora by any chance ?
How do I get a character to not be the main focus of an image? Trying to get some furry milfs on an alien planet. The milfs im the background and the general building/scenery in the foreground
Ive tried distant view and far away which did very little
no loras. im taking colored drawings and then doing img2img w them
controlnet openpose, the size difference betwen the "stick guy" and the size of the canvas is honored
maybe the amount of steps and/or upscaler ?
I never used SDXL before so I wanted to try today, I got a model and a VAE but it won't let me change the VAE, it keeps switching back to the one I used for SD 1.5.
No idea why
sdxl dosnt use vae as far as i know
how do I use it?
I tried to leave that VAE field on automatic but it takes ages to generate an image
i got 16GB VRAM, so usually it generates fast
This might be asking for the moon but is there a way to do this without controlnet? My limited vram cant do controlnet at the resolution I need
Is there an extension that is like photoshop generative fill?
you should be generating at 512x512 and then upscaling when the image is finished
Hires fix
512x512 produces awful images without atleast hires fix. 768x768 is way better with hires fix but that uses almost all my vram.
how long it takes for u to generate a 768x 768 with high res fix
Around 7 minutes with a 1660ti
thats just normal
Yup
mine takes 6mins for 1080p ( 544x960 x2 high res fix ) im also on 6g vram
u could try to use 1.5 high res fix then upscale on img2img ultimate scale sd
I can get a 1080p in roughly 10 minutes. I have to use tiles VAE tho.
Im surprised a 1660ti can do it at all tbh
hey at least you are not on amd 😬
True, my last card was an rtx 480 8gb, thought this one would be a slight downgrade on the gaming side but its good for AI
yea 16xx series is not that good for ai 10xx series is faster despite being an older generation dunno why
I'm eyeing a rtx a2000 12gb... I know it's older tech but I can run it in my SFF... any limitations to older cards that I might want to know?
Hello
I am new, why there is so many version for the model?
does everybody download model from civitai?
Hello, in which chat can I ask for help, advice on writing prom or other topics? (Nsfw)
hmm can't decide if I should get a (both are) samsung 1tb M2 PCIE4 for 100 euro or a 2tb M2 PCIE4 for 170 euro
leaning towards 2tb
win 11? System disk? @young ivy
already have one for system
I need room for games
and games are starting to require ssd or even m2s
in my case i cant use it as system disk, if you already have it then superb!
Hello, how can I add a bot to my discord server?
Samsung 990 Pro 2TB for 140€
Thats what I got for games and ai
Hey guys and gals.
Is there a way to load a Lora in A1 for the Hires Fix pass only?
Good morning, everyone! How are we all today?
When i will be able to pay to use it ?
does anyone know why i have to change the model before i start using a1111? if i don't change the model ill just get solid grey images and errors
what model gives you grey images with errors?
||∆||
any of them, if I just start a1111 and start going, without first changing the model
try update your extensions and restart if it helps? Not sure at all
na ive done those things
ive seen others w the problem
i had the problem on my last pc as well
a whole new build and fresh SD install and im still having it
hey hey!
Where can i ask a question about the software
@wild steppe sorry for interrupting but I got a question about sdxl
Welcome! Check out our starter channels such as a #1072229020520947753 or #1080946152318443610 and if you still have questions - #✨|sdxl would be the best place to ask
Ok I was just wondering what the terms and conditions for the software
Depends on where/how you're using it, it's best to refer to the applicable license
can someone help me
Anyone of you already tried a ghost mannequin effect?
in models/stable-diffusion
Where can i download musicgen to run on my own pc?
sorry for the interup, I'd type my prompt on channel bot, but its no response. And now I can't do any generate, what can I do?
I think that has to do with the bot that talks to the server which it is connected to and not Stable Diffusion but I might be gravely mistaken.
Thanks for answering! And It's about the bot. As I can't find any solution in FAQ Channel it's hard to find a place asking for helpT.T If I'm in the wrong channel please guide me!
Hi all. Can I add a bot to my discord server?
Hello, can I make a video with my photos?
Hey can someone recommend a good laptop that will run stable diffusion smoothly and fast enough?
Budget is 1k
Attention: If you or a loved one has been diagnosed with Mesothelioma you may to be entitled to financial compensation. Mesothelioma is a rare cancer linked to asbestos exposure. Exposure to asbestos in they Navy, shipyards, mills, heating, construction or the automotive industries may put you at risk.
Sorry, just ignore that, I didn't know it was fake
apparently google indexes bard conversation. You can view random ones just by searching for site:https://bard.google.com/share
needless to say, that's uh, not good
you are told the conversations are not secret, or im reading your coment wrong ?
most private google product
I can mix facial expressions by just adding more to the prompt? That what it feels like, not sure if its true.
Hey folks, updated both Automatic1111 and openOutpaint and now there is no more "Send to openOutpaint" buttons on the PNG tab (or anywhere else for that matter). Anyone else having this issue? Any suggestions on how to get the button back? There doesn't seem to be a way to load a PNG directly into openOutpaint that I can see...
are you not in incognito mode? I used to have issue i was unable to import images due this incognito mode.
so suggestion, try different browser if you have
I'm not in incognito mode, just using Chrome normally. I've never had this issue before with openOutpaint and been using it for months and months. But that's interesting to hear that incognito mode had that effect. Even though I'm not using it, I'll still try another browser and see what happens!
Opened it in MS edge but there was no difference
No worries and thank you for helping regardless (still haven't tried comfyui, should do that some day)
some nodes are realy interesting
Yeah I've heard good things about it!
Ok I resolved the issue. So... I updated automatic1111 like 2 days ago and removed git pull from the bat file because I want to be in control of when it updates. I didn't think that in the span of 2 days there would have been an update that broke it, so I didn't even try to update since 2 days ago. But I decided to try, put git pull back, it updates automatic1111 and voila! Now the button is BACK
So lesson learned: Always try the very very freshest update to see if that fixes things facepalms
😄 superb it works.
Sorry to bother you all with this. But perhaps if someone else runs into the same problem this might help them. At least when I searched I saw @undone garden had been asking the exact same thing in #🤝|tech-support . So Ciccio if you see this message, try updating everything.
Im having problem with using animated diff with comfyui, I am following this workflow:
https://civitai.com/articles/2314
after downloading all the modules and checkpoints im getting this error
Error occurred when executing KSampler: mat1 and mat2 shapes cannot be multiplied (1232x768 and 1024x320)
Could anyone give me any guidance?
Hi I am missing the SD VAE drop down on the top left corner for Stable diffusion. Is this an extension I have to dl? my sd version is 1.6.0
mat1 and mat2 seems you are using 1.5 model for sdxl thing or other way round. @waxen gull
sd vae is posible to download on civitai.com for example. For model you are using. 1,5 for 1,5 sdxl for sdxl.
In your link is a copy
How can I increase the token count of the prompt in the sdxl model to 90 tokens?
from diffusers import StableDiffusionXLPipeline
import torch
from PIL import Image
import openai
import re
import time
import os
prompt_list = []
for i, prompt in enumerate(prompts):
prompt_list.append(prompt)
print(prompt)
pipeline = StableDiffusionXLPipeline.from_pretrained(
"stabilityai/stable-diffusion-xl-base-1.0", torch_dtype=torch.float16, variant="fp16", use_safetensors=True
).to("cuda")
for i, prompt in enumerate(prompt_list):
generated_image = pipeline(prompt=prompt).images[0]
filename_prompt = prompt[:100]
output_path = f"{folder_path}/{filename_prompt}.jpg"
generated_image.save(output_path)
print(f"Generated image saved to {output_path}")
counter += 1
print(counter)```
my prompts looks like this:
A serene lakeside landscape at sunset, with a lone boat drifting on the water. It's a high-resolution photograph intended to feel natural, calming, and reflective. You can find it on Shutterstock or Getty Images, featuring extremely detailed and sharp focus. The image showcases warm, pastel colors and exudes an atmosphere of tranquility. The primary colors are rich golden tones with a hint of lavender. The image is accompanied by soft, warm sunlight casting long shadows.
Off to my third cancer surgery, third time it's back and this time it's not a certain outcome...Wish me luck in a way that won't jinx me
which nvidia driver version do you use ? 531.79 ? or latest ? 537.42 ?
Yo I need help
NansException: A tensor with all NaNs was produced in Unet. This could be either because there's not enough precision to represent the picture, or because your video card does not support half type. Try setting the "Upcast cross attention layer to float32" option in Settings > Stable Diffusion or using the --no-half commandline argument to fix this. Use --disable-nan-check commandline argument to disable this check.
Do u know how to solve this?
hey would need the full error of the cmd and your gpu name in #🤝|tech-support
I figured out why i was getting green highlights if anyone is interested
and why I had to change my model on UI start everytime to get an image
Still can't get my inpaintanything extension tab to show up. Anyone else had this issue?
Hello I have a question : image created with stable diffusion are 100% unique ? And can i use it for commercial use ?
a
hi guys
checking specs for a new pc, will these two give entirely different results ?
PNY GEFORCE RTX™ 4070Ti 12GB
PNY GeForce RTX™4080 16GB
Also is there a better or preffered company for AI stuff ?
dont go with 4070ti
what do you mean ?
Not go with 4070ti, it has issues, and it will make difference for future for sure 16GB is more futureproof
oh ok cool, good to know, any other models that could be better for the price ?
that's the specs they offered me for the 4080
ANTEC DF700 FLUX RGB
ASUS PRIME Z790-P WIFI Intel LGA1700 Z790 DDR5
INTEL I7-13700KF Tray No FAN/GPU 16Crs max
Corsair H150 RGB 360mm Liquid CPU Cooler
PNY GeForce RTX™4080 16GB VERTO™ Triple Fan
Corsair VENGEANCE 32GB 2x16GB DDR5 DRAM 5600MHz C36
Gigabyte AORUS Gen4 5000E SSD 1TB
Antec HCG GOLD 1000W Fully Modular 2X8Eps 8X8 PCI-E
around 3300$
where do you download the controlnet models?
guys, why is it that when I generate beautiful and high-quality art in Stable Diffusion, a new one appears on my half of the download lane, with defects and bugs? How to fix it?
if you have errors feel free to provide screenshots and info in #🤝|tech-support
setup is good but sadly 4080 is overpriced :/
I dont think the prices will go lower soon 😦
AI will keep it going up
true 🥲
then i'll probably go with this one lol

